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Open-source latent diffusion text-to-image model for high-quality, customizable image generation and editing.
Open-source latent diffusion text-to-image model for high-quality, customizable image generation and editing.
Stable Diffusion is a latent diffusion text-to-image model that converts text prompts (and images) into high-quality, photorealistic or stylized images. It conditions generation on text embeddings (commonly CLIP ViT-L/14) to guide image synthesis and supports variants such as Stable Diffusion XL (SDXL) and SDXL Turbo that prioritize improved composition, realism, and shorter prompt requirements. The model is distributed as open-source weights and is designed for fine-tuning, custom checkpoints, image-to-image transformations, inpainting/outpainting, and integration into hosted APIs and third-party tools. Its combination of open licensing, extensibility (fine-tuning, ControlNet, upscalers), and widespread integrations makes it a foundational model for creative workflows, production image pipelines, and research.